# Lightweight diffusion model

An anime-style LoRa adapter trained based on FLUX-1 development version, capable of generating high-quality anime-style images through trigger words

This is an OpenVINO-optimized midjourney-mini model for text-to-image generation tasks.

A LoRA fine-tuned model based on Stable Diffusion 3.5 Turbo, specifically designed for generating simple doodle-style images

A text-to-image generation model based on SDXL-Turbo, which combines LCM (Latent Consistency Models) and LoRA (Low-Rank Adaptation) technologies to achieve fast and high-quality image generation.

Lightweight UNet2D model trained on MNIST dataset for unconditional generation of handwritten digit images

A lightweight text-to-image generation model optimized through distillation based on Realistic_Vision_V4.0, achieving 80% faster speed than base SD1.5

Chilloutmix is a text-to-image generation model based on Stable Diffusion, suitable for generating high-quality artistic images.

A Stable Diffusion model fine-tuned using the DreamBooth method, specializing in text-to-image generation tasks.